This is an open-weight Stable Diffusion XL model specialized for inpainting. It allows users to generate or edit image content within masked regions of an input image using text prompts. The model is distributed via the Diffusers library on Hugging Face and can be run locally or via inference providers for image generation tasks.
In the Image generation space, Stable Diffusion Xl 1.0 Inpainting takes a focused approach. It focuses on generating or editing images by filling masked regions using text prompts. Stable Diffusion Xl 1.0 Inpainting is an open-source project aimed at developers. Stable Diffusion Xl 1.0 Inpainting is open source under the Open Source license. It ships for the web, the command line, and API.
Behind Stable Diffusion Xl 1.0 Inpainting is Hugging Face, and it first shipped in 2021. Development happens publicly on GitHub with 14k stars and 115 commits in the last 90 days. Among its 3 catalogued features are inpainting, text-to-Image, and Stable Diffusion XL. It exposes integrations via a public API.
Summary written by a language model from the project’s public pages.
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