This is an open-weight Stable Diffusion XL model specialized for inpainting. It allows users to generate or edit image content within masked regions of an input image using text prompts. The model is distributed via the Diffusers library on Hugging Face and can be run locally or via inference providers for image generation tasks.
Stable Diffusion Xl 1.0 Inpainting is an Image generation product. It focuses on generating or editing images by filling masked regions using text prompts. It is built as an open-source project for developers. Stable Diffusion Xl 1.0 Inpainting is open source under the Open Source license. Stable Diffusion Xl 1.0 Inpainting is available on the web, the command line, and API.
Hugging Face builds and maintains Stable Diffusion Xl 1.0 Inpainting, and the product first shipped in 2021. Development happens publicly on GitHub with 14k stars and 115 commits in the last 90 days. Key capabilities include inpainting, text-to-Image, and Stable Diffusion XL. It exposes integrations via a public API.
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